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Fooocus Image Prompt

All text to image models

Default parameters with automated optimizations and quality improvements.

InferenceFrom 1 credit per run

About Fooocus Image Prompt

Default parameters with automated optimizations and quality improvements.

stylized

Pricing

A run with the default settings costs 1 credits. The exact cost updates in the playground as you change settings, and you're only charged when you press Run. Failed runs are refunded automatically.

Inputs

prompttextareadefault: ""

The prompt to use for generating the image. Be as descriptive as possible for best results.

negative_prompttextareadefault: ""

The negative prompt to use. Use it to address details that you don't want in the image. This could be colors, objects, scenery and even the small details (e.g. moustache, blurry, low resolution).

stylesstring-listdefault: ["Fooocus V2","Fooocus Sharp","Fooocus Enhance"]

The style to use.

performanceselectdefault: "Extreme Speed"

You can choose Speed or Quality

Options: Speed, Quality, Extreme Speed, Lightning

guidance_scalenumberdefault: 4

The CFG (Classifier Free Guidance) scale is a measure of how close you want the model to stick to your prompt when looking for a related image to show you.

sharpnessnumberdefault: 2

The sharpness of the generated image. Use it to control how sharp the generated image should be. Higher value means image and texture are sharper.

aspect_ratiotextdefault: "1024x1024"

The size of the generated image. You can choose between some presets or custom height and width that **must be multiples of 8**.

num_imagesintegerdefault: 1

Number of images to generate in one request

lorasjsondefault: [{"path":"https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/resolve/main/sd_xl_offset_example-lora_1.0.safetensors","scale":0.1}]

The LoRAs to use for the image generation. You can use up to 5 LoRAs and they will be merged together to generate the final image.

refiner_modelselectdefault: "None"

Refiner (SDXL or SD 1.5)

Options: None, realisticVisionV60B1_v51VAE.safetensors

refiner_switchnumberdefault: 0.8

Use 0.4 for SD1.5 realistic models; 0.667 for SD1.5 anime models 0.8 for XL-refiners; or any value for switching two SDXL models.

output_formatselectdefault: "jpeg"

The format of the generated image.

Options: png, jpeg, webp

sync_modebooleandefault: false

If `True`, the media will be returned as a data URI and the output data won't be available in the request history.

seedinteger

The same seed and the same prompt given to the same version of Stable Diffusion will output the same image every time.

image_prompt_1jsonrequired
image_prompt_2json
image_prompt_3json
image_prompt_4json
inpaint_image_urlfile

The image to use as a reference for inpainting.

mask_image_urlfile

The image to use as a mask for the generated image.

inpaint_modeselectdefault: "Inpaint or Outpaint (default)"

The mode to use for inpainting.

Options: Inpaint or Outpaint (default), Improve Detail (face, hand, eyes, etc.), Modify Content (add objects, change background, etc.)

inpaint_additional_prompttextareadefault: ""

Describe what you want to inpaint.

outpaint_selectionsstring-listdefault: []

The directions to outpaint.

mixing_image_prompt_and_inpaintbooleandefault: false

Mixing Image Prompt and Inpaint

uov_image_urlfile

The image to upscale or vary.

uov_methodselectdefault: "Disabled"

The method to use for upscaling or varying.

Options: Disabled, Vary (Subtle), Vary (Strong), Upscale (1.5x), Upscale (2x), Upscale (Fast 2x)

mixing_image_prompt_and_vary_upscalebooleandefault: false

Mixing Image Prompt and Vary/Upscale

enable_safety_checkerbooleandefault: true

If set to false, the safety checker will be disabled.

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